An Introduction to Diffusion Models and Stable Diffusion

An Introduction to Diffusion Models and Stable Diffusion

Introduction

Imagine a world where creativity transcends the limitations of brushes, clay, and canvas. It was in 2022, at Colorado’s State Fair art competition, that a groundbreaking entry defied the conventional boundaries of artistic creation. Jason M. Allen’s masterpiece, “Théâtre D’opéra Spatial” won the first prize and defied convention. Not through traditional means, but with the aid of an AI program called Midjourney, which uses a diffusion model to generate images. By turning a text prompt into a hyper-realistic image, Allen’s creation not only captivated the audience and judges but also set off a fierce backlash from artists who accused him of, essentially, cheating.

“Théâtre D’opéra Spatial” entry for the Colorado State Fair.

However, the rise of Midjourney and other AI advancements merely scratches the surface of what is possible with diffusion models. These generative models have become a force to be reckoned with, attracting attention and pushing the boundaries of image synthesis previously ruled by Generative Adversarial Networks (GANs) [1]. In fact, diffusion models are being said to be beating GANs on image synthesis [2].

In this article, we explore the theoretical foundations of diffusion models, uncovering their inner workings and understanding their fundamental components and remarkable effectiveness. Along the way, we’ll shine a spotlight on one of the most popular families of diffusion models: Stable Diffusion.

Join us as we uncover the secrets behind diffusion models’ success and how they are revolutionizing image generation. By the end, you’ll have a profound understanding of their transformative potential, inspiring the realm of AI-driven creativity.

Diffusion Models

Diffusion Models are generative models that learn from the data during training and generate similar examples based on what they have learned. These models draw inspiration from non-equilibrium thermodynamics and have achieved state-of-the-art quality in generating various forms of data. Some examples include generating high-quality images and even audio (e.g. Audio Diffusion Models [3]). If you are interested in how diffusion models can be applied in audio settings, check our blog on voice cloning

In a nutshell, Diffusion Models work by corrupting training data through the addition of Gaussian noise (called forward diffusion process), and then learning how to recover the original information by reversing this noising process step by step (called reverse diffusion process). Once trained, these models can generate new data by sampling random Gaussian noise and passing it through the learned denoising process.

Diffusion Models go beyond just creating high-quality images. They have gained popularity by addressing the well-known challenges associated with adversarial training in GANs. Diffusion Models offer advantages in terms of training stabilityefficiencyscalability, and parallelization.

In the following sections, we’ll dive deeper into the fine details of Diffusion Models. We’ll explore the forward diffusion process, the reverse diffusion process, and an overview of the steps involved in the training process. Also, we’ll also gain an intuition on the calculation of the loss function. By examining these components, we’ll acquire a comprehensive understanding of how Diffusion Models function and how they achieve their impressive results.

Forward Diffusion Process

The forward diffusion process consists of gradually adding Gaussian noise to an input image step by step, for a total of T steps. At step 0 we have the original image, at step 1 a very slightly corrupted image which gets even more corrupted step by step until the whole information of the original image is lost.

Forward diffusion process [4]

To formalize this process, we can view it as a fixed Markov chain with T steps, where the image at timestep t maps to its subsequent state at timestep t+1. As such, each step depends only on the previous one, allowing us to derive a closed-form formula to obtain the corrupted image at any desired timestep, bypassing the need for iterative computation.

Forward diffusion formulation [4]

Consequently, this closed-form formula enables direct sampling of xₜ at any timestep, significantly accelerating the forward diffusion process.

Schedulers

In addition, the noise addition at each step follows a deliberate pattern. A Scheduler determines the amount of noise to be added. In the original Denoising Diffusion Probabilistic Models (DDPM) paper [4], the authors define a linear schedule ranging βₜ from 0.0001 at timestep 0 to 0.02 at timeste

### Stable Diffusion 模型概述 Stable Diffusion 是一种先进的图像生成模型,其核心在于潜在扩散模型(LDM),该模型能够在较低维度的潜在空间中进行操作[^1]。这种特性使得模型能够高效地处理复杂的图像生成任务,在保持高质量的同时显著降低了计算资源的需求。 #### 架构组成 Stable Diffusion 的架构主要由三个部分构成: - **CLIP Model**: 负责将文本描述转换为向量形式,并将其映射至与图像相同的嵌入空间,以便于后续处理阶段中的匹配和对比[^2]。 - **UNet**: 作为去噪网络的一部分,负责接收来自 VAE 编码器压缩后的特征图谱,并逐步去除其中的人工噪声,最终恢复出清晰的目标图像[^3]。 - **VAE (Variational Autoencoder)**: 实现了从原始像素级表示到紧凑潜在表征间的双向变换功能,既可完成降维编码也可逆向解码重建原貌。 ### 应用实例 为了便于实际部署和调用,开发者通常会构建一条完整的流水线来简化交互流程。下面是一个简单的 Python 示例代码片段展示了如何利用 Hugging Face 提供的 `transformers` 库快速启动一个基于 Stable Diffusion 的文本转图片服务[^4]: ```python from transformers import pipeline, Tasks import torch pipe = pipeline( task=Tasks.text_to_image_synthesis, model='multi-modal_chinese_stable_diffusion_v1.0', torch_dtype=torch.float16 if torch.cuda.is_available() else torch.float32 ) result = pipe("A beautiful sunset over the ocean") ``` 此段脚本定义了一个名为 `pipe` 的对象,代表了一条集成了预处理、推理以及后处理环节在内的自动化生产线。当给定一段文字说明时,程序将会依据内置算法自动生成相应的视觉效果。
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